La ecuación de la recta que contiene el punto (x, y) = (0, 6) y es paralela a la recta 2 · x + 3 · y + 4 = 0 es igual a 2 · x + 3 · y - 18 = 0.
¿Cómo determinar la ecuación de una recta?
En esta pregunta tenemos que derivar una ecuación de la recta a partir de una recta y un punto. Dos rectas bajo forma implícita son paralelas si y solo si comparten los mismos coeficientes dependientes, entonces tenemos que el coeficiente independiente de la recta es:
2 · 0 + 3 · 6 + k = 0
18 + k = 0
k = - 18
Por tanto, la ecuación de la recta que contiene el punto (x, y) = (0, 6) y es paralela a la recta 2 · x + 3 · y + 4 = 0 es igual a 2 · x + 3 · y - 18 = 0.
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Find the angle of elevation of the sun from the ground when a tree that is 18 ft tall casts a shadow 25 ft long. Round to the nearest degree.
Answer:
36°
Step-by-step explanation:
Let [tex]\theta[/tex] be the angle of elevation. The side opposite of [tex]\theta[/tex] will be the height of the tree, which is 18ft, and the side adjacent to [tex]\theta[/tex] will be the length of the shadow, which is 25ft. Because these two lengths are known, then we should use the tangent ratio to determine the measure of the angle of elevation:
[tex]\displaystyle \tan\theta=\frac{\text{Opposite}}{\text{Adjacent}}=\frac{18}{25}\biggr\\\\\\\theta=\tan^{-1}\biggr(\frac{18}{25}\biggr)\approx36^\circ[/tex]
Therefore, the angle of elevation is about 36°.
Solve: 3x - 3 = x + 1
Hello !
Answer:
[tex]\Large\boxed{ \sf x = 2}[/tex]
Step-by-step explanation:
Let's solve the following equation by isolating x.
[tex] \sf3x - 3 = x + 1[/tex]
First, add 3 to both sides :
[tex] \sf3x - 3 + 3 = x + 1 + 3[/tex]
[tex] \sf3x = x + 4[/tex]
Now let's substract x from both sides :
[tex] \sf3x - x = 4[/tex]
[tex] \sf2x = 4[/tex]
Finally, let's divide both sides by 2 :
[tex] \sf \frac{2x}{2} = \frac{4}{2} [/tex]
[tex] \boxed{ \sf x = 2}[/tex]
Have a nice day ;)
Students from Logan, Kennedy, Newark Memorial, and Hayward High have Debate teams in the finals. List the possible ways the four schools can place 1st, 2nd, 3rd, and 4th.
There are 24 possible ways the four schools can place 1st, 2nd, 3rd, and 4th in the finals.
We have,
To determine the possible ways the four schools can place 1st, 2nd, 3rd, and 4th in the finals, we can use the concept of permutations.
A permutation is an arrangement of objects in a specific order. In this case, we want to find the permutations of the four schools.
The number of permutations can be determined by multiplying the number of choices for each position.
Since there are four schools, there are four choices for the 1st position, three choices for the 2nd position, two choices for the 3rd position, and one choice for the 4th position.
The total number of permutations is given by:
= 4 × 3 × 2 × 1
= 24
Therefore,
There are 24 possible ways the four schools can place 1st, 2nd, 3rd, and 4th in the finals.
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Henry needs to give informal proof of the formula for the circumference of a circle.
He first constructs a circle, with center O, and labels a point on the circle as P.
He draws a radius from O to P.
He then uses point P as the center to construct a new circle.
He draws two line segments, each formed by joining point O with the points of intersection of the two circles.
Which of these is a plausible next step in Henry's proof process?
Construct another circle with a doubled radius.
Construct a rectangle that circumscribes the original circle.
Construct an octagon that circumscribes the original circle.
Construct a hexagon inscribed in the original circle
The circumference of a circle is given by the following formula:
C = 2πr
Where C is the circumference and r is the radius of the circle.
Henry has constructed a circle, with center O, and labeled a point on the circle as P.
He has drawn a radius from O to P and used point P as the center to construct a new circle.
He has drawn two line segments, each formed by joining point O with the points of intersection of the two circles.
A plausible next step in Henry's proof process is to construct a rectangle that circumscribes the original circle.
Circumscribing a circle means creating a geometric figure that encloses the given circle but does not have any overlapping points.
A circle circumscribed inside a rectangle is shown in the figure below:
A circle can also be circumscribed by polygons, such as an equilateral triangle, a square, a regular hexagon, and so on.
In this case, the polygon is drawn so that each vertex of the polygon touches the circle.
The circumference of a circle is given by the following formula:
C = 2πr
Where C is the circumference and r is the radius of the circle.
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y′′ 2y′ y = f(t), y(0) = 0, y′(0) = 1, where f(t) = { 0 0 ≤t < 3 2 3 ≤t <10, f(t) = 0,t > 10
The solution to the given differential equation is y(t) = te^t + 1/2 for 3 ≤ t < 10, and y(t) = c1 e^t + c2 t e^t for t < 3 and t ≥ 10.Note that the constants c1 and c2 in the last expression can be determined using the continuity of y(t) and y′(t) at t = 3 and t = 10.
To solve the given differential equation, we first find the general solution to the homogeneous equation y′′ - 2y′ + y = 0. The characteristic equation is r^2 - 2r + 1 = 0, which has a double root of r = 1. Therefore, the general solution to the homogeneous equation is y_h(t) = (c1 + c2 t)e^t.
Next, we find a particular solution to the non-homogeneous equation. Since f(t) is piecewise defined, we consider two cases:
Case 1: 0 ≤ t < 3. In this case, f(t) = 0, so the non-homogeneous equation becomes y′′ - 2y′ + y = 0. We already have the general solution to this equation, so the particular solution is y_p(t) = 0.
Case 2: 3 ≤ t < 10. In this case, f(t) = 2, so the non-homogeneous equation becomes y′′ - 2y′ + y = 2. We try a particular solution of the form y_p(t) = At + B. Substituting this into the equation gives A = 0 and B = 1/2. Therefore, the particular solution in this case is y_p(t) = 1/2.
Case 3: t ≥ 10. In this case, f(t) = 0, so the non-homogeneous equation becomes y′′ - 2y′ + y = 0. We already have the general solution to this equation, so the particular solution is y_p(t) = 0.
The general solution to the non-homogeneous equation is y(t) = y_h(t) + y_p(t) = (c1 + c2 t)e^t + 1/2 for 3 ≤ t < 10, and y(t) = (c1 + c2 t)e^t for t < 3 and t ≥ 10.Using the initial conditions, we have y(0) = 0, which implies that c1 = 0. Also, y′(0) = 1, which implies that c2 = 1.
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prove that, for all integers m and n, 4 | (m2 n 2 ) if and only if m and n are even.
We have proved both implications, we can conclude that 4 divides (m^2 * n^2) if and only if m and n are both even.
How to prove m and n are even?To prove that 4 divides (m^2 * n^2) if and only if m and n are even, we need to prove two implications:
If 4 divides (m^2 * n^2), then m and n are even.
If m and n are even, then 4 divides (m^2 * n^2).
Let's start with the first implication:
If 4 divides (m^2 * n^2), then m and n are even.
We can prove this by contrapositive. Assume that m and n are not both even, which means that at least one of them is odd. Without loss of generality, let's assume that m is odd. Then m can be written as m = 2k + 1, where k is an integer. Substituting this into the expression for m^2 * n^2, we get:
m^2 * n^2 = (2k + 1)^2 * n^2
= 4k^2 * n^2 + 4kn^2 + n^2
Note that the first two terms in this expression are both divisible by 4, but the last term (n^2) is not necessarily divisible by 4, since n could be odd. Therefore, m^2 * n^2 is not divisible by 4 if m and n are not both even. This proves the contrapositive, and hence the first implication.
Now, let's move on to the second implication:
If m and n are even, then 4 divides (m^2 * n^2).
We can prove this directly. Since m and n are even, we can write them as m = 2k and n = 2j, where k and j are integers. Substituting these into the expression for m^2 * n^2, we get:
m^2 * n^2 = (2k)^2 * (2j)^2
= 4k^2 * 4j^2
= 16(k^2 * j^2)
Since k and j are integers, k^2 * j^2 is also an integer, and hence 16(k^2 * j^2) is divisible by 4. Therefore, m^2 * n^2 is divisible by 4 if m and n are both even. This proves the second implication.
Since we have proved both implications, we can conclude that 4 divides (m^2 * n^2) if and only if m and n are both even.
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find the radius of convergence, r, of the series. [infinity] (8x 5)n n2 n = 1 r = 1 8 find the interval, i, of convergence of the series. (enter your answer using interval notation.) i =
The given series is:
∑(n=1 to ∞) (8x^5)^n/n^2
We can use the ratio test to determine the radius of convergence:
lim (n→∞) |(8x^5)^(n+1)/(n+1)^2| / |(8x^5)^n/n^2|
= lim (n→∞) (8x^5)/(n+1)^2 * n^2/(8x^5)
= lim (n→∞) n^2/(n+1)^2
= 1
The radius of convergence is:
r = 1/8
To find the interval of convergence, we need to test the endpoints x = -r and x = r:
When x = -r = -1/8:
∑(n=1 to ∞) (8(-1/8)^5)^n/n^2 = ∑(n=1 to ∞) (-1)^n/n^2
This is an alternating series with decreasing terms, so we can use the alternating series test to show that it converges. Therefore, the series converges when x = -1/8.
When x = r = 1/8:
∑(n=1 to ∞) (8(1/8)^5)^n/n^2 = ∑(n=1 to ∞) 1/n^2
This is a convergent p-series with p = 2, so it converges by the p-series test. Therefore, the series converges when x = 1/8.
The interval of convergence is therefore:
i = [-1/8, 1/8]
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Because a p-value of zero, while theoretically possible is effectively impossible, a p-value of .00000 is written as
A. < .01
B. 0.01
C. ~ .00000
D. Approximately .00000
E. All of the above
F. none of the above
The p-value of zero, while theoretically possible is effectively impossible, a p-value of .00000 is written as is
option A, "< .01".
It is important to first understand what a p-value represents. A p-value is a statistical measure that indicates the likelihood of obtaining the observed results of a study or experiment by chance, assuming that there is no true effect or difference between groups.
In hypothesis testing, a p-value of less than .05 (or .01, depending on the level of significance chosen) is typically considered to be statistically significant, indicating that the observed results are unlikely to be due to chance alone.
However, a p-value of exactly zero is not possible, as it would mean that the observed results are absolutely certain and could not have occurred by chance. Therefore, a p-value of .00000 (or any other extremely small value) is typically reported as "< .01" or something similar, indicating that the p-value is less than the chosen level of significance (in this case, .01).
Therefore, option A, "< .01", is the most accurate way to represent a p-value of .00000. The other options are either not precise enough (B and D), or incorrect (C and F).
A p-value of exactly zero is impossible, and a p-value of .00000 (or any other extremely small value) is typically reported as "< .01" to indicate that it is less than the chosen level of significance.
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Write an exponential function in the form y=ab^xy=ab
x
that goes through points (0, 7)(0,7) and (5, 1701)(5,1701)
To write an exponential function in the form y = ab^x that passes through the given points (0, 7) and (5, 1701), we can use these points to find the values of a and b.
Let's start by substituting the coordinates of the first point (0, 7) into the equation:
7 = ab^0
7 = a
So we have determined that a = 7.
Now, let's substitute the coordinates of the second point (5, 1701) into the equation:
1701 = 7b^5
To isolate b, we can divide both sides of the equation by 7:
1701/7 = b^5
Now, we can simplify the left side of the equation:
243 = b^5
Taking the fifth root of both sides, we find:
b = 3
Therefore, we have determined that a = 7 and b = 3.
Putting it all together, the exponential function that goes through the given points is:
y = 7 * 3^x
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Let f(x) = tan x. a) show that f is 1-1 and differentiable on (-pi/2, pi/2), hence has a differentiable inverse. b) Let g denote the inverse. Use the inverse function theorem to find g'(y) for any real y.
The result g'(y) = cos^2 g(y) for any real y.
To show that f(x) = tan x is 1-1 and differentiable on (-pi/2, pi/2), we can use the fact that the derivative of tan x is sec^2 x, which is continuous and positive on (-pi/2, pi/2).
This means that f(x) is increasing and never constant on this interval, thus satisfying the 1-1 condition. Furthermore, since sec^2 x is continuous on this interval, f(x) is also differentiable.
To find the inverse function g'(y), we can use the inverse function theorem, which states that if f is differentiable and 1-1 in an open interval containing x and if f'(x) is not equal to 0, then its inverse function g is differentiable at y = f(x) and g'(y) = 1/f'(x). Applying this theorem to f(x) = tan x, we have:
f'(x) = sec^2 x
f'(g(y)) = sec^2 g(y)
g'(y) = 1/f'(g(y)) = 1/sec^2 g(y) = cos^2 g(y)
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g'(y) = cos^2(g(y)) for any real y. This formula gives us the derivative of the inverse function g(x) of f(x) = tan(x).
a) To show that f(x) = tan(x) is one-to-one (1-1) on the interval (-π/2, π/2), we need to demonstrate that for any two distinct values of x in the interval, their corresponding function values are also distinct.
Let x1 and x2 be two distinct values in (-π/2, π/2), such that x1 ≠ x2. We then have:
f(x1) = tan(x1) and f(x2) = tan(x2)
To prove that f is 1-1, we need to show that if f(x1) = f(x2), then x1 = x2. Taking the contrapositive, if x1 ≠ x2, then f(x1) ≠ f(x2).
Assume x1 ≠ x2. We know that the tangent function has a period of π, so the values of tan(x) repeat after every π units. However, since x1 and x2 are both in the interval (-π/2, π/2), their corresponding tangent values will be distinct. Therefore, f(x1) ≠ f(x2), and we have shown that f is 1-1 on (-π/2, π/2).
To show that f is differentiable on (-π/2, π/2), we can demonstrate that the derivative of f(x) = tan(x) exists and is continuous on the interval. The derivative of tan(x) is sec^2(x), which is defined and continuous on (-π/2, π/2). Hence, f(x) = tan(x) is differentiable on (-π/2, π/2).
b) Since f(x) = tan(x) is 1-1 and differentiable on (-π/2, π/2), it has a differentiable inverse denoted as g(x).
According to the inverse function theorem, if f is differentiable and 1-1 on an interval I, and if f'(x) ≠ 0 for all x in I, then g'(y) = 1 / f'(g(y)).
In this case, f(x) = tan(x), which has a derivative of f'(x) = sec^2(x). Since f'(x) ≠ 0 for all x in (-π/2, π/2), we can use the inverse function theorem to find g'(y) for any real y.
Using the formula g'(y) = 1 / f'(g(y)), we substitute f(x) = tan(x) and solve for g'(y):
g'(y) = 1 / f'(g(y))
g'(y) = 1 / sec^2(g(y))
g'(y) = cos^2(g(y))
Therefore, g'(y) = cos^2(g(y)) for any real y. This formula gives us the derivative of the inverse function g(x) of f(x) = tan(x).
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You are planning to make an open rectangular box from a 10 inch by 19 inch piece of cardboard by cutting congruent squares from thr corners and folding up the sides.
What are the dimensions of the box of largest volume you can make this way, and what is its volume?
Length = 19 - 2x ≈ 11.334 inches
Width = 10 - 2x ≈ 2.334 inches
Height = x ≈ 3.833 inches
V ≈ 167.386 cubic inches
Let x be the side length of each square cut from the corners of the cardboard. Then the length, width, and height of the resulting box will be:
Length = 19 - 2x
Width = 10 - 2x
Height = x
The volume of the box is given by:
V = length × width × height
V = (19 - 2x) × (10 - 2x) × x
Expanding the product and simplifying, we get:
V = 4x^3 - 58x^2 + 190x
To find the value of x that maximizes the volume, we can take the derivative of V with respect to x and set it equal to zero:
dV/dx = 12x^2 - 116x + 190 = 0
Solving for x using the quadratic formula, we get:
x = (116 ± sqrt(116^2 - 4×12×190)) / (2×12) ≈ 3.833 or 7.833
Since x must be less than 5 (half the width of the cardboard), the only valid solution is x ≈ 3.833.
Therefore, the dimensions of the box of largest volume are:
Length = 19 - 2x ≈ 11.334 inches
Width = 10 - 2x ≈ 2.334 inches
Height = x ≈ 3.833 inches
And its volume is:
V ≈ 167.386 cubic inches
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A 10m ladder is leaning against house the base of the ladder is pulled away from the houseat a rate of 0. 25m/sec how fast is the top of the ladder moving down the wall when the base is8m from the house?
The top of the ladder is moving down the wall at a rate of approximately 0.67 m/sec when the base is 8m from the house.
The height of the ladder is 10m.
The rate of the ladder base moving away from the wall is 0.25m/sec.
The distance between the ladder base and the wall is 8m.
We need to find how fast the top of the ladder is moving down the wall when the base is 8m from the house.
Given that the rate of the ladder base moving away from the wall is 0.25m/sec, we can find the rate at which the top of the ladder is moving down the wall by using related rates theorem.
Let's call the distance between the top of the ladder and the ground "y" and the distance between the bottom of the ladder and the wall "x".
We can use the Pythagorean Theorem to relate x and y:y^2 + x^2 = 10^2.
Differentiating both sides of the equation with respect to time, we get:2yy' + 2xx' = 0
Rearranging the equation, we get: y' = -(xx')/y.
Plugging in the given values, we get: y' = -8(0.25)/sqrt(10^2 - 8^2)≈ -0.67 m/sec.
Therefore, the top of the ladder is moving down the wall at a rate of approximately 0.67 m/sec when the base is 8m from the house.
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7. In AABC, AC||DE. In ABCG, CG||EF. Prove that: AD:DB = GF:FB G F E D B
To prove that AD:DB = GF:FB, we will use the properties of parallel lines and their transversals.
Given:
AC || DE (Line AC is parallel to line DE)
CG || EF (Line CG is parallel to line EF)
We can start by applying the property of parallel lines and their transversals to triangle ABC and triangle EDC:
By the Intercept theorem, we have:
AD/DB = CE/ED ...(1)
Now, let's apply the property of parallel lines and their transversals to triangle BCG and triangle FED:
By the Intercept theorem, we have:
GF/FB = DE/EC ...(2)
Since AC || DE and CG || EF, we know that CE = AC and DE = CG. Therefore, we can substitute these values into equations (1) and (2):
From equation (1):
AD/DB = AC/CG
From equation (2):
GF/FB = CG/AC
Since AC = CE and CG = DE, we can rewrite these equations as:
AD/DB = CE/DE
GF/FB = DE/CE
Since DE = CE, we can conclude that:
AD/DB = GF/FB
Therefore, we have proved that AD:DB = GF:FB using the properties of parallel lines and their transversals.
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The plants in Tara's garden have a 6-foot x 10-foot area in which to grow. The garden is bordered by a brick walkway of width w.
Part A: Write two equivalent expressions to describe the perimeter of Tara's garden, including the walkway.
Part B: How can you check to see if your two expressions from Part A are equivalent?
Part C: What is the total perimeter of Tara's garden including the walkway if the walkway is 2.5ft wide?
The total perimeter of the garden is 42ft if the walkway is 2.5ft wide.
Part A:Two equivalent expressions to describe the perimeter of Tara's garden including the walkway are:
2(6 + w) + 2(10 + w) = 24 + 4w, where w is the width of the walkway.
The 2(6 + w) accounts for the two lengths of the rectangle, and 2(10 + w) accounts for the two widths of the rectangle. Simplify the expression to 4w + 24 to give the total perimeter of the garden. The other expression is:
20 + 2w + 2w + 12 = 2w + 32
Part B:To check the equivalence of the two expressions from Part A, we could simplify both expressions, as shown below.2(6 + w) + 2(10 + w) = 24 + 4w.
Simplifying the expression will yield:2(6 + w) + 2(10 + w)
= 2(6) + 2(10) + 4w2(6 + w) + 2(10 + w)
= 32 + 4w2(6 + w) + 2(10 + w)
= 4(w + 8)
Similarly, we can simplify 20 + 2w + 2w + 12 = 2w + 32, which yields:20 + 2w + 2w + 12 = 4w + 32
Part C:If the walkway is 2.5ft wide, the total perimeter of Tara's garden, including the walkway, is:
2(6 + 2.5) + 2(10 + 2.5)
= 2(8.5) + 2(12.5)
= 17 + 25
= 42ft.
We can find two equivalent expressions to describe the perimeter of Tara's garden, including the walkway. We can use the expression 2(6 + w) + 2(10 + w) and simplify it to 4w + 24.
The other expression can be obtained by adding the length of all four sides of the garden. We can check the equivalence of both expressions by simplifying each expression and verifying if they are equal.
We can calculate the total perimeter of Tara's garden, including the walkway, by using the formula 2(6 + 2.5) + 2(10 + 2.5), which gives us 42ft as the answer.
Thus, the conclusion is that the total perimeter of the garden is 42ft if the walkway is 2.5ft wide.
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Determine the convergence or divergence of the sequence with the given nth term. If the sequence converges, find its limit. (If the quantity diverges, enter DIVERGES.)
an = ln(n9) / 2n
To determine the convergence or divergence of the sequence with the given nth term:
The given nth term is: an = ln(n^9) / (2n)
As n approaches infinity, we can analyze the behavior of the sequence:
Taking the limit as n approaches infinity:
lim (n → ∞) ln(n^9) / (2n)
Using the properties of logarithms, we can rewrite the expression as:
lim (n → ∞) 9ln(n) / (2n)
Applying L'Hôpital's rule:
By differentiating the numerator and denominator with respect to n, we get:
lim (n → ∞) (9/n) / 2
Simplifying further:
lim (n → ∞) 9 / (2n)
As n approaches infinity, the term (2n) in the denominator grows indefinitely, causing the entire expression to converge to zero.
Therefore, the given sequence converges, and its limit is 0.
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Determine whether the geometric series is convergent or divergent. 10 - 6 + 18/5 - 54/25 + . . .a. convergentb. divergent
After applying the ratio test to the given geometric series, the answer is option a: the series is convergent.
Is the given geometric series convergent or divergent?The given series is: 10 - 6 + 18/5 - 54/25 + ...
To determine whether this series is convergent or divergent, we can use the ratio test.
The ratio test states that a series of the form ∑aₙ is convergent if the limit of the absolute value of the ratio of successive terms is less than 1, and divergent if the limit is greater than 1. If the limit is equal to 1, then the ratio test is inconclusive.
So, let's apply the ratio test to our series:
|ax₊₁ / ax| = |(18/5) * (-25/54)| = 15/20.24 ≈ 0.74
As the limit of the absolute value of the ratio of successive terms is less than 1, we can conclude that the series is convergent.
Therefore, the answer is (a) convergent.
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standard error is same as a. standard deviation of the sampling distribution b. difference between two means c. variance of the sampling distribution d. variance
The answer is that the standard error is the standard deviation of the sampling distribution. This means that it measures the amount of variability or spread in the means of multiple samples drawn from the same population.
To understand the concept of standard error, it is important to distinguish it from the standard deviation, which measures the amount of variability or spread in a single sample. The standard error, on the other hand, reflects the precision of the sample mean as an estimate of the population mean. It takes into account the fact that different samples will produce different means due to chance variation.
More specifically, the standard error is calculated by dividing the standard deviation of the population by the square root of the sample size. This formula reflects the fact that larger sample sizes tend to produce more precise estimates of the population mean, while smaller sample sizes are more likely to have greater sampling error or deviation from the true mean.
The standard error is used in many statistical analyses, particularly in hypothesis testing and constructing confidence intervals. For example, if we want to determine whether a sample mean is significantly different from a hypothesized population mean, we would calculate the standard error and use it to compute a t-value or z-value. This value would then be compared to a critical value to determine the statistical significance of the difference. Similarly, in constructing a confidence interval, we use the standard error to estimate the range of values that are likely to contain the true population mean with a certain level of confidence.
The standard error is the standard deviation of the sampling distribution, and it reflects the precision of the sample mean as an estimate of the population mean. It is calculated by dividing the standard deviation of the population by the square root of the sample size, and it is used in many statistical analyses to test hypotheses and construct confidence intervals.
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use the discriminant to determine whether the equation of the given conic represents an ellipse, a parabola, or a hyperbola. −6x2 4xy 12y2−9x 2y−8=0
The given equation represents an ellipse, since Δ = 304 is greater than zero.
The given equation, −6x^2 + 4xy + 12y^2 − 9x − 2y − 8 = 0, represents a second-degree equation involving both x and y. To determine the type of conic, we can analyze the discriminant. The discriminant is calculated as Δ = B^2 − 4AC, where A, B, and C are the coefficients of the x^2, xy, and y^2 terms, respectively.
In this case, A = -6, B = 4, and C = 12. Substituting these values into the discriminant formula, we get Δ = (4)^2 - 4(-6)(12) = 16 + 288 = 304.
By examining the value of the discriminant, we can classify the conic as follows:
- If Δ > 0, the conic is an ellipse.
- If Δ = 0, the conic is a parabola.
- If Δ < 0, the conic is a hyperbola.
Since Δ = 304, which is greater than zero, the given equation represents an ellipse.
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a sample has a mean of m = 86. if one new person is added to the sample, and σx is unchanged, what effect will the addition have on the sample mean?
As σx (standard deviation) remains unchanged, the value of x alone cannot determine the effect on sample mean. It depends on the value of x relative to the values in original sample and sample size.
If one new person is added to the sample and the standard deviation (σx) remains unchanged, the effect on the sample mean (m) can be determined as follows:
Let's denote the original sample size as n and the sum of the sample values as Σx.
Original sample mean:
m = Σx / n
After adding the new person, the new sample size becomes n + 1, and the sum of the sample values becomes Σx + x_new (x_new represents the value of the new person).
New sample mean:
m' = (Σx + x_new) / (n + 1)
To analyze the effect, we can express the difference in means:
Δm = m' - m = ((Σx + x_new) / (n + 1)) - (Σx / n)
Simplifying this expression, we get:
Δm = (x_new - (Σx / n)) / (n + 1)
Therefore, the effect of adding the new person on the sample mean (m) is determined by the difference between the value of the new person (x_new) and the original mean (Σx / n), divided by the increased sample size (n + 1).
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Using the same context as the previous problem - A toy race car is racing on a circular track and the car is 4 feet from the center of the racetrack. After only traveling around 80% of the track, the motor in the car stopped working and the toy race car was stuck. a. How far along the track (in feet) did the toy race car travel before stopping? b. How many radians did the toy race car sweep out from its starting position to when it stopped working? c. How far is the toy race car to the right of the center of the track (in feet) when it traveled 80% of the track? d. If the toy race car travels an additional 2radians from where it stopped working on the track how far will the toy race car be to the right of the center of the track? e. If the toy race car travels an additional 4T radians from where it stopped working on the track how far will the toy race car be to the right of the center of the track?
a) The toy race car traveled 20.106 feet before stopping.
b) The toy race car swept out approximately 1.6π radians from its starting position to when it stopped working.
c) The toy race car is 4 feet to the left of the center of the track when it traveled 80% of the track.
d) The toy race car will be approximately 1.236 feet to the right of the center of the track.
e) The toy race car will be at x = 4 cos(4T + 1.6π) + 2.55 to the right of the center of the track.
a. To find how far along the track the toy race car traveled before stopping, we can simply multiply the circumference of the circular track by 0.8, since the car traveled 80% of the track before stopping.
Circumference = 2πr
= 2π(4) (since the car is 4 feet from the center of the track)
= 8π feet
Distance traveled = 0.8 × 8π
= 6.4π feet
= 20.106 feet (rounded to three decimal places)
Therefore, the toy race car traveled 20.106 feet before stopping.
b. To find how many radians the toy race car swept out from its starting position to when it stopped working, we can use the formula:
θ = s/r
where θ is the angle in radians, s is the distance traveled along the arc, and r is the radius of the circle.
We know that the distance traveled along the arc is 0.8 times the circumference of the circle, which we calculated to be 8π feet. The radius of the circle is 4 feet. Therefore:
θ = (0.8 × 8π) / 4
= 1.6π radians
= 5.026 radians (rounded to three decimal places)
Therefore, the toy race car swept out 5.026 radians from its starting position to when it stopped working.
c. To find how far the toy race car is to the right of the center of the track when it traveled 80% of the track, we need to find the horizontal displacement of the toy race car at that point. Since the toy race car is traveling on a circular track, we can use trigonometry to find its horizontal displacement.
The distance traveled by the toy race car along the track is 80% of the circumference of the circle, which is:
circumference = 2πr = 2π(4) = 8π feet
distance traveled = 0.8 × 8π = 6.4π feet
This distance corresponds to an angle of:
angle = distance traveled / radius = 6.4π / 4 = 1.6π radians
Using this angle, we can find the horizontal displacement using cosine:
cos(1.6π) = -1
Therefore, the toy race car is 4 feet to the left of the center of the track when it traveled 80% of the track.
d. To find how far the toy race car will be to the right of the center of the track if it travels an additional 2 radians from where it stopped working, we can use the same trigonometric approach as in part c. We know that the radius is 4 feet and the toy race car will sweep out an additional angle of 2 radians, so its horizontal displacement will be:
cos(1.6π + 2) = -cos(0.4π) = -0.309
Therefore, the toy race car will be approximately 1.236 feet to the right of the center of the track.
e. If the toy race car travels an additional 4T radians from where it stopped working on the track, we can use the same approach as in part d. The position of the toy race car is given by:
x = r cos(θ) + d
where θ = 4T and d is the distance from the center of the track (found in part c). Plugging in the values, we get:
x = 4 cos(4T + 1.6π) + 2.55
Note that the value of x will depend on the value of T, which is not given in the problem.
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Generate a number that has a digit in the tenths place that is 100 times smaller than the 8 in the hundreds place. 184. 36
A number that has a digit in the tenths place that is 100 times smaller than the 8 in the hundreds place is 184.36.
Let's break down the given number, 184.36. The digit in the hundreds place is 8, which is 100 times larger than the digit in the tenths place.
In the decimal system, each place value to the right is 10 times smaller than the place value to its immediate left. Therefore, the digit in the tenths place is 100 times smaller than the digit in the hundreds place. In this case, the tenths place has the digit 3, which is indeed 100 times smaller than 8.
So, by considering the value of each digit in the number, we find that 184.36 satisfies the condition of having a digit in the tenths place that is 100 times smaller than the 8 in the hundreds place.
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Consider two independent continuous random variables X1, X2 each uniformly distributed over [0, 2]. Let Y = max (X1, X2), i.e., the maximum of these two random variables. Also, let Fy (y) be the cumulative distribution function (CDF) of Y. Find Fy (y) where y = 0.72.
The CDF of Y evaluated at y = 0.72 is 0.1296.
Since X1 and X2 are independent and uniformly distributed over [0, 2], their joint density function is:
f(x1, x2) = 1/4, for 0 ≤ x1 ≤ 2 and 0 ≤ x2 ≤ 2
To find the CDF of Y, we can use the fact that:
Fy(y) = P(Y ≤ y) = P(max(X1, X2) ≤ y)
This event can be split into two cases:
X1 and X2 are both less than or equal to y:
In this case, Y will be less than or equal to y.
The probability of this occurring can be calculated using the joint density function:
P(X1 ≤ y, X2 ≤ y) = ∫0y ∫0y f(x1, x2) dx1 dx2
= ∫0y ∫0y 1/4 dx1 dx2
[tex]= (y/2)^2[/tex]
[tex]= y^2/4[/tex]
One of X1 or X2 is greater than y:
In this case, Y will be equal to the maximum of X1 and X2.
The probability of this occurring can be calculated as the complement of the probability that both X1 and X2 are less than or equal to y:
P(X1 > y or X2 > y) = 1 - P(X1 ≤ y, X2 ≤ y)
[tex]= 1 - y^2/4[/tex]
Therefore, the CDF of Y is:
Fy(y) = P(Y ≤ y) = P(max(X1, X2) ≤ y)
= P(X1 ≤ y, X2 ≤ y) + P(X1 > y or X2 > y)
[tex]= y^2/4 + 1 - y^2/4[/tex]
= 1, for y ≥ 2
[tex]= y^2/4,[/tex]for 0 ≤ y ≤ 2
To find Fy(0.72), we simply substitute y = 0.72 into the expression for Fy(y):
[tex]Fy(0.72) = (0.72)^2/4 = 0.1296[/tex]
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Since X1 and X2 are uniformly distributed over [0, 2], their probability density functions (PDFs) are:
fX1(x) = fX2(x) = 1/2, for 0 <= x <= 2
To find the CDF of Y = max(X1, X2), we need to consider two cases:
1. If y <= 0, then Fy(y) = P(Y <= y) = 0
2. If 0 < y <= 2, then Fy(y) = P(Y <= y) = P(max(X1, X2) <= y)
We can find this probability by considering the complementary event, i.e., the probability that both X1 and X2 are less than or equal to y. Since X1 and X2 are independent, this probability is:
P(X1 <= y, X2 <= y) = P(X1 <= y) * P(X2 <= y) = (y/2) * (y/2) = y^2/4
Therefore, the CDF of Y is:
Fy(y) = P(Y <= y) =
0, y <= 0
y^2/4, 0 < y <= 2
1, y > 2
To find Fy(0.72), we substitute y = 0.72 into the CDF:
Fy(0.72) = 0.72^2/4 = 0.1296
Therefore, the value of Fy(y) at y = 0.72 is 0.1296.
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Jonah's monthly salary is R. 13200. If 12% is deducted for tax,1% for UIF and 2% for pension, how much does jonah receive each month after deductions?
Jonah will receive R 11 320 each month after all the deductions. Jonah's monthly salary is R. 13200. If 12% is deducted for tax,1% for UIF and 2% for pension, the amount that Jonah receives each month after the deductions will be: Firstly, let's calculate the amount that Jonah will be taxed.
Jonah's monthly salary is R. 13200. If 12% is deducted for tax,1% for UIF and 2% for pension, the amount that Jonah receives each month after the deductions will be: Firstly, let's calculate the amount that Jonah will be taxed. For this, we will multiply his salary by the percentage that will be deducted for tax: 12/100 x 13200 = R 1584
Next, we will calculate the amount that Jonah will pay for UIF. For this, we will multiply his salary by the percentage that will be deducted for UIF: 1/100 x 13200 = R 132
Finally, we will calculate the amount that Jonah will pay for pension. For this, we will multiply his salary by the percentage that will be deducted for pension: 2/100 x 13200 = R 264
Total amount that will be deducted = R 1980
Amount that Jonah will receive after deductions = R 13200 - R 1980 = R 11 320
Therefore, Jonah will receive R 11 320 each month after all the deductions. This question deals with calculating the monthly salary of Jonah after the deductions.
The problem stated that Jonah's monthly salary is R. 13200. It was further stated that 12% of his salary is deducted for tax, 1% for UIF and 2% for pension. From the given information, we have to calculate the amount that Jonah receives each month after the deductions.To solve the problem, we started by calculating the amount that will be deducted for tax. For this, we multiplied Jonah's salary by the percentage that will be deducted for tax i.e 12/100. The product of these two values came out to be R 1584.Then, we calculated the amount that Jonah will pay for UIF. For this, we multiplied his salary by the percentage that will be deducted for UIF i.e 1/100. The product of these two values came out to be R 132.
Finally, we calculated the amount that Jonah will pay for pension. For this, we multiplied his salary by the percentage that will be deducted for pension i.e 2/100. The product of these two values came out to be R 264.The total amount that will be deducted is the sum of the values that we calculated above. Therefore, the total amount that will be deducted is R 1980.To find out the amount that Jonah will receive each month after the deductions, we subtracted the total amount of the deductions from his monthly salary. The result of this calculation came out to be R 11 320. Therefore, Jonah will receive R 11 320 each month after all the deductions.
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The AO, of Adequate intake of water, for pregnant women is a mean of 3L/d, liters per day. Sample data n=200, x=2. 5, s=1. The sample data appear to come from a normally distributed population with a 0=1. 2
The sample mean is 2.5 liters per day, and the sample standard deviation is 1 liter. The population mean is given as 3 liters per day. It appears that the sample data come from a normally distributed population.
The sample data provides information about the daily water intake of pregnant women. The sample size is 200, and the sample mean is 2.5 liters per day, with a sample standard deviation of 1 liter. The population mean, or Adequate Intake (AI), for pregnant women is given as 3 liters per day.
To determine if the sample data come from a normally distributed population, additional information is required. In this case, the population standard deviation is not provided, but the population mean is given as 3 liters per day.
If the sample data come from a normally distributed population, we can use statistical tests such as the t-test or confidence intervals to make inferences about the population mean. However, without additional information or assumptions, we cannot conclusively determine if the sample data come from a normally distributed population.
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2. LetA=\begin{bmatrix} a &b \\ c & d \end{bmatrix}(a) Prove that A is diagonalizable if (a-d)2 + 4bc > 0 and is not diagonalizable if (a-d)2 + 4bc < 0.(b) Find two examples to demonstrate that if (a-d)2 + 4bc = 0, then A may or may not be diagonalizble.
We can find the eigenvalues of [tex]$A$[/tex] using the characteristic equation:
[tex]$$\det(A-\lambda I) = \begin{vmatrix} a-\lambda & b \\ c & d-\lambda \end{vmatrix} = (a-\lambda)(d-\lambda) - bc = \lambda^2 - (a+d)\lambda + (ad-bc)$$[/tex]
The discriminant of this quadratic equation is:
[tex]$$(a+d)^2 - 4(ad-bc) = (a-d)^2 + 4bc$$[/tex]
Therefore, [tex]$A$[/tex] is diagonalizable if and only if [tex]$(a-d)^2 + 4bc > 0$[/tex].
If [tex]$(a-d)^2 + 4bc > 0$[/tex], then the discriminant is positive, and the characteristic equation has two distinct real eigenvalues. Since [tex]$A$[/tex] has two linearly independent eigenvectors, it is diagonalizable.
If [tex]$(a-d)^2 + 4bc < 0$[/tex], then the discriminant is negative, and the characteristic equation has two complex conjugate eigenvalues. In this case, [tex]$A$[/tex] does not have two linearly independent eigenvectors, and so it is not diagonalizable.
(b) If [tex]$(a-d)^2 + 4bc = 0$[/tex], then the discriminant of the characteristic equation is zero, and the eigenvalues are equal. We can find two examples to demonstrate that [tex]$A$[/tex] may or may not be diagonalizable in this case.
Example 1: Consider the matrix [tex]$A = \begin{bmatrix} 1 & 2 \\ 2 & 4 \end{bmatrix}$[/tex]. We have [tex]$(a-d)^2 + 4bc = (1-4)^2 + 4(2)(2) = 0$[/tex], so the eigenvalues of [tex]$A$[/tex] are both [tex]$\lambda = 2$[/tex]. The eigenvectors are [tex]$\begin{bmatrix} 1 \\ 1 \end{bmatrix}$[/tex] and [tex]$\begin{bmatrix} -2 \\ 1 \end{bmatrix}$[/tex], respectively. Since these eigenvectors are linearly independent, [tex]$A$[/tex] is diagonalizable.
Example 2: Consider the matrix [tex]$A = \begin{bmatrix} 1 & 1 \\ -1 & -1 \end{bmatrix}$[/tex]. We have [tex]$(a-d)^2 + 4bc = (1+1)^2 + 4(-1)(-1) = 0$[/tex], so the eigenvalues of[tex]$A$[/tex] are both [tex]$\lambda = 0$[/tex]. The eigenvector is[tex]$\begin{bmatrix} 1 \\ -1 \end{bmatrix}$[/tex], which is the only eigenvector of [tex]A$. Since $A$[/tex] has only one linearly independent eigenvector, it is not diagonalizable.
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use the ratio test to determine whether the series is convergent or divergent. [infinity]
∑ (−1)^n−1 (9^n / 5^n n^3)
n = 1 identify an
The ratio test is a useful tool for determining the Confluence or divergence of a series. It involves taking the limit of the absolute value of the rate of the( n 1) th term to the utmost term as n approaches perpetuity. Depending on the value of this limit, we can determine whether the series converges or diverges.
The rate test is a important tool used to determine the confluence or divergence of a series. It involves taking the limit of the absolute value of the rate of the( n 1) th term to the utmost term as n approaches infinity. However, also the series converges absolutely, If this limit is lower than 1. still, also the series diverges, If the limit is lesser than 1. still, also the test is inconclusive and another test must be used, If the limit equals 1. Using the rate test, we can determine the confluence or divergence of a given series. For illustration, if we've a series( perpetuity) n = 1 of a_n, also we can apply the rate test by taking the limit as n approaches perpetuity of| a,( n 1)/a_n|.
still, also the series converges absolutely, If this limit is lower than 1. still, also the series diverges, If the limit is lesser than 1. still, also we can not determine the confluence or divergence of the series using the rate test alone, If the limit equals 1. the rate test is a useful tool for determining the confluence or divergence of a series. It involves taking the limit of the absolute value of the rate of the( n 1) th term to the utmost term as n approaches perpetuity. Depending on the value of this limit, we can determine whether the series converges or diverges.
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The limit of the ratio test is (9/5), which is less than 1. Therefore, by the ratio test, the series converges.
We can use the ratio test to determine the convergence of the series:
|(-1)^(n+1) (9^(n+1) / 5^(n+1) (n+1)^3)| / |(-1)^(n) (9^n / 5^n n^3)|
= (9/5) * (n^3/(n+1)^3)
Taking the limit as n approaches infinity:
lim (n^3/(n+1)^3) = lim (1/(1+1/n))^3 = 1
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-4d^-3 simplify the expression so all exponents are positive
To simplify the expression and make all exponents positive, we can use the rule that says that a negative exponent is the same as the reciprocal of the corresponding positive exponent. In other words,
a^(-n) = 1/(a^n)
Using this rule, we can rewrite the given expression as:
-4d^-3 = -4/(d^3)
Therefore, the simplified expression with all exponents positive is -4/(d^3).
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Show that (A) if A and B are Hermitian, then AB is not Hermitian unless A and B commute (B) a product of unitary matrices is unitary
A) If A and B are Hermitian, then AB is not Hermitian unless A and B commute.
B) A product of unitary matrices is unitary.
A) Proof:
Let A and B be Hermitian matrices. Then, A and B are defined as A* = A and B* = B.
We know that the product of two Hermitian matrices is not necessarily Hermitian, unless they commute. This means that AB ≠ BA.
Thus, if A and B do not commute, then AB is not Hermitian.
B) Proof:
Let U and V be two unitary matrices. We know that unitary matrices are defined as U×U=I and V×V=I, where I denotes an identity matrix.
Then, we can write the product of U and V as UV = U*V*V*U.
Since U* and V* are both unitary matrices, the product UV is unitary as U*V*V*U
= (U*V*)(V*U)
= I.
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(A) If A and B are Hermitian matrices that do not commute, AB is not Hermitian.
(B) The product of two unitary matrices, UV, is unitary.
Let's begin with statement (A):
(A) If A and B are Hermitian, then AB is not Hermitian unless A and B commute.
To prove this statement, we will use the fact that for a matrix to be Hermitian, it must satisfy A = A^H, where A^H denotes the conjugate transpose of A.
Assume that A and B are Hermitian matrices. We want to show that if A and B do not commute, then AB is not Hermitian.
Suppose A and B do not commute, i.e., AB ≠ BA.
Now let's consider the product AB:
(AB)^H = B^H A^H [Taking the conjugate transpose of AB]
Since A and B are Hermitian, we have A = A^H and B = B^H. Substituting these in, we get:
(AB)^H = B A
If AB is Hermitian, then we should have (AB)^H = AB. However, in general, B A ≠ AB unless A and B commute.
Therefore, if A and B are Hermitian matrices that do not commute, AB is not Hermitian.
Now let's move on to statement (B):
(B) A product of unitary matrices is unitary.
To prove this statement, we need to show that the product of two unitary matrices is also unitary.
Let U and V be unitary matrices. We want to show that UV is unitary.
To prove this, we need to demonstrate two conditions:
1. (UV)(UV)^H = I [The product UV is normal]
2. (UV)^H(UV) = I [The product UV is also self-adjoint]
Let's analyze the two conditions:
1. (UV)(UV)^H = UVV^HU^H = U(VV^H)U^H = UU^H = I
Since U and V are unitary matrices, UU^H = VV^H = I. Therefore, (UV)(UV)^H = I.
2. (UV)^H(UV) = V^HU^HU(V^H)^H = V^HVU^HU = V^HV = I
Similarly, since U and V are unitary matrices, V^HV = U^HU = I. Therefore, (UV)^H(UV) = I.
Thus, both conditions are satisfied, and we conclude that the product of two unitary matrices, UV, is unitary.
In summary:
(A) If A and B are Hermitian matrices that do not commute, AB is not Hermitian.
(B) The product of two unitary matrices, UV, is unitary.
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TRUE/FALSE. Ap-value is the highest level (of significance) at which the observed value of the test statistic is insignificant.
The statement is true because the p-value represents the highest level of significance at which the observed value of the test statistic is considered insignificant.
When conducting hypothesis testing, the p-value is calculated as the probability of obtaining a test statistic as extreme as, or more extreme than, the observed value, assuming the null hypothesis is true. It is compared to the predetermined significance level (alpha) chosen by the researcher.
If the p-value is greater than the chosen significance level (alpha), it indicates that the observed value of the test statistic is not statistically significant. In this case, we fail to reject the null hypothesis, as the evidence does not provide sufficient support to reject it.
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Solve the differential equation. Mention the method you used.
(t2+1)dydt=yt−y(2+1)
After solving the given differential equation, the resultant answer is: y = C(t^2+1)^{1/2} (t-3)/(2t)
To solve the differential equation (t2+1)dydt=yt−y(2+1), we need to first separate the variables y and t. We do this by dividing both sides of the equation by (yt-y(2+1))/(t^2+1) to get:
dy/(y - (2+1)/t) = dt/(t^2+1)
Now, we use the method of partial fractions to simplify the left-hand side of the equation. Let's write:
y/(y - (2+1)/t) = A + Bt/(t^2+1)
where A and B are constants that we need to find. Multiplying both sides by the denominator (y - (2+1)/t)(t^2+1), we get:
y = A(y^2+1) + Bt(y-(2+1)/t)
Expanding and collecting like terms, we get:
y^2(A-B/t) + Bt^2y - B(t+2) = 0
Since this equation must hold for all y and t, we can equate coefficients of y^2, y, and the constant term to get three equations:
A - B/t = 0
Bt^2 = 1
-B(t+2) = 0
From the second equation, we get B = 1/t^2. Substituting this into the first equation, we get A = 1/t. Finally, the third equation tells us that B = 0 if t = -2. Thus, the partial fraction decomposition is:
y/(y - (2+1)/) = 1/t + t/(t^2+1)
Substituting this into the separated differential equation, we get:
(1/t + t/(t^2+1)) dy = dt/(t^2+1)
Integrating both sides, we get:
ln|y - (2+1)/t| + 1/2 ln(t^2+1) = ln|C| + arctan(t)
where C is the constant of integration. Solving for y, we get:
y = C(t^2+1)^{1/2} (t-3)/(2t)
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